Stable diffusion xl prompt guide. Keep reading to start creating.

HiRes: 4xNMKD-Siax_200k with 3 Steps and 0. Minimalist background, 50mm lens for focused facial expression. CFG 16 – 20: Not generally recommended unless the prompt is well-detailed. Might affect coherence and quality. A collection of what Stable Diffusion imagines these artists' styles look like. Negative prompting influences the generation process by acting as a high-dimension anchor, which Mar 16, 2024 · For local AI image generation, it’s hard to beat Stable Diffusion. Faster examples with accelerated inference. Parameters . Relevant: Use relevant keywords and phrases that are related to the subject and Cool Stable Diffusion XL prompts. Switch between documentation themes. 0) Image folder: <path to your image folder> Output folder: <path to the Nov 30, 2023 · prompt #10: enchanting and colorful wonderland with fantastical creatures, whimsical landscapes, and playful adventures. I’ve covered vector art prompts, pencil illustration prompts, 3D illustration prompts, cartoon prompts, caricature prompts, fantasy illustration prompts, retro illustration prompts, and my favorite, isometric illustration prompts in this Oct 29, 2023 · This Stable Diffusion XL or SDXL prompt guide aims to provide a comprehensive understanding of various aspects related to the SDXL model from prompt anatomy to styles. Oct 24, 2023 · There is an option in Stable Diffusion that enables you to select this aspect ratio. Sep 16, 2023 · A negative prompt is a way to use Stable Diffusion in a way that allows the user to specify what he doesn’t want to see, without any extra input. Stable Diffusion XL (SDXL) 1. You can use the free AI image generator on Stable Diffusion Online or search over 9 million Stable Diffusion prompts on Prompt Database. Stable Diffusion XL can pass a different prompt for each of the text encoders it was trained on as shown below. Examples of prompts for the Stable Diffusion process. In this list, you’ll find various styles you can try with SDXL models. New stable diffusion finetune ( Stable unCLIP 2. The following prompts are mostly collected from different discord servers, websites, fabricated and then modified Mar 19, 2024 · There are two main ways to make videos with Stable Diffusion: (1) from a text prompt and (2) from another video. You can use them as is or as a starting point to create your own. The prompt affects the output for a trivial reason. Let’s make some images Aug 4, 2023 · Undoubtedly, the emergence of Stable Diffusion XL has marked a milestone in the history of natural language processing and image generation, taking us a step closer to something that already scares… Sep 16, 2023 · In this comprehensive guide, we’ll walk you through setting up the software, using the color sketch tool, and leveraging Img2Img to turn amateur sketches into professional artwork. Copy it to your favorite word processor, and apply it the same way as before, by pasting it into the Prompt field and clicking the blue arrow button under Generate. First, we establish a simple image generation workflow using the 🧨 Diffusers library by 🤗 HuggingFace. Fooocus. 0 model without any LORA models. - keep a file of prompt ideas that you have copied and try them out. 0 specifically for pixel art. co/ponyxl_loras_n_stuff was used to source as well as the purplesmart. It is a parameter that tells the Stable Diffusion model what not to include in the generated image. She wears a medieval dress. 5 Upscale. On the checkpoint tab in the top-left, select the new “sd_xl_base” checkpoint/model. Please post them in plain text over here with their accompanying pictures - that would be much more useful and reach more people. 5 and 768×768 for SD 2. You have probably seen one of them on social media. Aug 19, 2023 · Ce guide a pour vocation de vous aider à maîtriser l'interface graphique d'AUTOMATIC1111. Aug 6, 2023 · Learn how to use SDXL v1. Compared to Stable Diffusion V1 and V2, Stable Diffusion XL has made the following optimizations: Improvements have been made to the U-Net, VAE, and CLIP Text Encoder components of Stable Diffusion. We advise a fragmented Prompt such as: “ 1girl, schoolgirl, white uniform “ rather than a full sentence like: “ a schoolgirl in white uniform “ Even though they give very similar results with very small prompts, long sentence-type prompts are prone to be partially dismissed or get interrupted by unintended filler words. to get started. 9: Good luck, and always be testing! Stable Diffusion XL (SDXL) is an open-source diffusion model, the long waited upgrade to Stable Diffusion v2. You have to be specific. Written by: Milan K. Prompts On the first tab, txt2img, you'll be making most of your images. 0 model is built on innovation, comprising a 3. Aug 5, 2023 · Stable Diffusion XL can produce images at a resolution of up to 1024×1024 pixels, compared to 512×512 for SD 1. Il est conçu pour servir de tutoriel, avec de nombreux exemples pour illustrer l’utilité ou le fonctionnement d’un paramètre. I mentioned the Stable Diffusion XL model a few times in this guide. 2. What is Img2Img in Stable Diffusion. or. Adding noise in a specific order governed by Gaussian distribution concepts is Feb 22, 2024 · Stable Diffusion XL 1. Prompt examples : Prompt: cartoon character of a person with a hoodie , in style of cytus and deemo, ork, gold chains, realistic anime cat, dripping black goo, lineage revolution style, thug life, cute anthropomorphic bunny, balrog, arknights, aliased, very buff, black and red and yellow paint, painting illustration collage style, character Oct 31, 2023 · A negative prompt for SDXL is like giving it a description of what you don’t want to see in the picture. Example 1: An image of a beautiful sunset at the beach might be tagged by CLIP as: ocean, beach, sunset, waves, sand, palm trees, golden sky. Prompt: A beautiful ((Ukrainian Girl)) with very long straight hair, full lips, a gentle look, and very light white skin. Fooocus is a rethinking of Stable Diffusion and Midjourney’s designs: Learned from Stable Diffusion, the software is offline, open source, and free. No. prompt (str or List[str], optional) — The prompt or prompts to guide the image generation. Diffusion models are taught by introducing additional pixels called noise into the image data. Conclusion Oct 28, 2022 · Prompt Engineering Intro This is an example of a prompt and all the parameters Prompt: Funko pop superman figurine, made of plastic, product studio shot, on a white background, diffused lighting, centered. Step 2. You can harness these prompts as provided or adapt them as inspiration for your personalized phantom. Keep reading to start creating. Stable Diffusion XL output image can be improved by making use of a refiner as shown Aug 10, 2023 · SDXL 是 Stable Diffusion XL 的簡稱，顧名思義它的模型更肥一些，但相對的繪圖的能力也更好。 Stable Diffusion XL - SDXL 1. Enter your prompt in the top one and your negative prompt in the bottom one. Best. Stable Diffusion Web UI is a browser interface based on the Gradio library for Stable Diffusion. It's designed for designers, artists, and creatives who need quick and easy image creation. Summary. In this blog, we'll explore captivating prompts, each unlocking the potential of the SDXL 1. photograph should not be used with van Gogh. 0 can be accessed and used at no cost. The structure of the prompt A Detailed Stable Diffusion Pose Prompt Guide with 15 Examples. So, that’s our list of the best landscape prompts for Stable Diffusion. fashion editorial, a female model with blonde hair, wearing a colorful dress. Feb 12, 2024 · Here is our list of the best portrait prompts for Stable Diffusion: S. Related: Best Stable Diffusion Anime Prompts. 1, Hugging Face) at 768x768 resolution, based on SD2. Tags vs Sentence. 1st row prompt: <framing>, BREAK a gorgeous woman with short dark hair, new york city, <framing>, sunset Stable Diffusion models. 05 times. Now Stable Diffusion returns all grey cats. This applies to anything you want Stable Diffusion to produce, including landscapes. 6B parameter The following prompts are supposed to give an easier entry into getting good results in using Stable Diffusion. We deployed our model here; You can even run your trained model on segmind instead of the API; Fine-tuned model deployed as an API on Segmind. links and info on use https:// rentry. Table of Contents. These tips and prompting styles will work with any model that directly uses pony diffusion v6 xl, like autismix pony for example. Clipdrop is a fantastic way to test out Stable Diffusion XL, for free! They do offer a pro version for faster image generation, but even on the free version the images don’t take very long at all to be generated. Mar 29, 2024 · The fundamental usage of Stable Diffusion models is in the text-to-image format. 0 Model - Stable Diffusion XL Stable Diffusion XL or SDXL is the latest image generation model that is tailored towards more photorealistic outputs… Start by loading up your Stable Diffusion interface (for AUTOMATIC1111, this is “user-web-ui. Running Stable Diffusion Locally. It starts by creating functions to tokenize the prompts to calculate the prompt embeddings, and to compute the image embeddings with the VAE. The jump from 768×768 in SD 2. Feb 29, 2024 · Andrew. Oct 22, 2023 · For the purpose of the demonstration we will use the same prompt and use the XYZ Plot option in Stable Diffusion WebUI to plot both the camera framing and angle. blurry, noisy, deformed, flat, low contrast, unrealistic, oversaturated, underexposed. . I’ve categorized the prompts into different categories since digital illustrations have various styles and forms. Add a Comment. Stable Diffusion image 1 using 3D rendering. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. Stable Diffusion XL enables us to create gorgeous images with shorter descriptive prompts, as well as generate words within images. It uses a larger base model, and an additional refiner model to increase the quality of the base model’s output. L'ajout de quelques prompts négatifs peut améliorer considérablement l'esthétique des images générées. I'll be using that style for every prompt featured in this article. Photo of a man with a mustache and a suit, plain background, portrait style. This guide will focus on the code that is unique to the SDXL training script. Snowy Monster. We can even pass different parts of the same prompt to the text encoders. Jul 2, 2023 · A good Stable Diffusion prompt should be: Clear and specific: Describe the subject and scene in detail to help the AI model generate accurate images. We’re on a journey to advance and democratize artificial intelligence through open source and open science. Jan 25, 2024 · Poised, mid-dance movement, creating a sense of fluidity and grace. Be detailed and specific when describing the subject. This model allows for image variations and mixing operations as described in Hierarchical Text-Conditional Image Generation with CLIP Latents, and, thanks to its modularity, can be combined with other models such as KARLO. Jan 12, 2023 · Stable Diffusion is a text-to-image model that uses a frozen CLIP ViT-L/14 text encoder. We present SDXL, a latent diffusion model for text-to-image synthesis. Sign Up. Prompt: Monochromatic portrait of a model with dramatic makeup, chiaroscuro lighting for a classic art feel. However, you can use Lora models to generate images with a specific style, object, or setting. Embarking on a journey with Stable Diffusion prompts necessitates an exploratory approach towards crafting veritably articulate and distinctly specified prompts. It can create images in variety of aspect ratios without any problems. It provides a user-friendly way to interact with Stable Diffusion, an open-source text-to-image generation model. Very specific. NFTs seemingly came out of nowhere and skyrocketed in popularity. The more vivid your mental image, the more detailed your prompt can be. That's why I wrote two prompts to go along with that specific style. 500. In Stable Diffusion v1, most users do that with multiple custom models. Stable Diffusion is a Dec 18, 2023 · Enter any default prompt and click Import, the model will be imported in 2-3 minutes. The “ Icons. ago. ; prompt_2 (str or List[str], optional) — The prompt or prompts to be sent to the tokenizer_2 and text_encoder_2. The model distinguishes between concepts like “The Red Square” (a renowned location) and a “red square” (a geometric shape). We also demonstrate the use of the Weights & Biases integration for Diffusers to May 16, 2024 · Recommended Setting Hyper Version (VAE is baked in): Res: 832*1216 (Any SDXL Res will work fine) Sampler: DPM++ SDE Karras. 0 as its base checkpoint, which is the cutting edge Stable Diffusion model at the time of writing, it is a versatile LoRA that can do both app No Account Required! Stable Diffusion Online is a free Artificial Intelligence image generator that efficiently creates high-quality images from simple text prompts. 0 API Guide. Find recommended settings, samplers, VAE, and 33 examples of prompt templates for various genres. Now, let me start showing you the examples I wrote for you today. Dec 21, 2023 · Step-by-Step Guide for Stable Diffusion XL Setup. It looks like this. Simple prompts can already lead to good outcomes, but sometimes it's in the details on what makes an image believable. This compendium, which distills insights gleaned from a multitude of experiments and the collective wisdom of fellow Stable Diffusion aficionados, endeavors to be a Nov 20, 2023 · Best Stable Diffusion XL Photorealistic Prompts. Stable Diffusion Portrait Prompts. ← ControlNet with Stable Diffusion 3 ControlNet-XS →. . Use multiple brackets () to increase its strength and [] to reduce. Since it is open source and anyone who has 5GB of GPU VRAM can download it (and Emad Mostaque Oct 30, 2023 · Today, after Stable Diffusion XL is out, the model understands prompts much better. 0 is the latest model in the Stable Diffusion family of text-to-image models from Stability AI. 0 to 1024×1024 in SDXL represents a significant increase in the number of pixels – nearly Jan 22, 2024 · Approach 1: SD XL Generated App Icons. Begin by confirming the installation of Python, as Python is the backbone of this project then, install or update Pip, the essential package manager for Python, to ensure all dependencies for Stable Diffusion XL can be managed effectively. The latest version of this model is Stable Diffusion XL, which has a larger UNet backbone network and can generate even higher quality images. You can keep adding descriptions of what you want, including accessorizing the cats in the pictures. Stable Diffusion XL is the latest iteration of the popular text-to-image generation model, offering impressive results. Artist Inspired Styles. We would like to show you a description here but the site won’t allow us. Stable Diffusion XL. Aug 30, 2022 · In this guide, you will learn how to write prompts by example. I think using a video to present prompts is just about the worst medium you could use. Following this, an optional refiner model can be applied, which is responsible for adding more intricate details to the image. Use Detailed Subjects and Scenes to Make Your Stable Diffusion Prompts More Specific. Sep 23, 2023 · tilt-shift photo of {prompt} . It reminds me the moment when Stability AI Jan 29, 2024 · Food Photography. But crafting a great prompt takes time and effort and sometimes it may not be enough because language and words can be imprecise. - the UI, model, image dimensions, seed and other factors determine if your image is going to look like their image. Img2Img is a cutting-edge technique that generates new images from an input image and a corresponding text prompt. The architecture of the SDXL 1. I ultimately decided that I was going to use the 'Fantasy Art' preset style in Stable Diffusion XL. To access SDXL using Clipdrop, follow the steps below: Navigate to the official Stable Diffusion XL page on Clipdrop. 5B parameter base model and a 6. It is created by Stability AI. bat”). safetensors (you can also use stable-diffusion-xl-base-1. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. Jul 6, 2024 · You should see two nodes labeled CLIP Text Encode (Prompt). One of the standout features of this model is its ability to create prompts based on a keyword. In order to increase emphasis on a word or phrase, add a + or number between 1. In Stable Diffusion, however, triple parentheses amplify a word’s Nov 23, 2023 · I really like the fact that there is a preset style in Stable Diffusion XL 1. Good balance between creativity and guided generation. With dynamic batching enabled, concurrent requests, and info-level logging, the example prints out additional information about the number of prompts included in each request to the TensorRT Prompt techniques. Fooocus is an image generating software (based on Gradio ). Use "Cute grey cats" as your prompt instead. 3D rendering. I batch up the prompt and leave my PC for 30-40mins while it churns out the images. In this post, you will learn how it works, how to use it, and some common use cases. Asian Zen Garden. Depending on the hardware available to you, this can be very computationally intensive and it may not run on a Collaborate on models, datasets and Spaces. Published on: November 12, 2023. 0. Stable Diffusion image 2 using 3D rendering. Jun 22, 2023 · Want to learn prompting techniques within Stable Diffusion to produce amazing results from your ideas? Well, look no further than this short, straight to the Rule 2. In order to reduce emphasis on a word or phrase, add a - or number between 0 and 0. In all seriousness, models can be trained using different data sets to excel in specific types of tasks. Jan 27, 2024 · Is Lora required to create landscape art in Stable Diffusion? Lora models aren’t required to create landscape art. If you're unsure how to select this style, I'll show you in the image below. Through posture, artists can profoundly depict a character's personality, emotions, and inner world, providing viewers with a deeper understanding of the artwork. Not Found. N'hésitez pas à le bookmarquer pour le consulter également comme un manuel de référence. The second way is to stylize a video using Stable Diffusion. The more weight you add, the greater the risk of lower quality there is. ). Golden Retriever. Setup Python and Pip. selective focus, miniature effect, blurred background, highly detailed, vibrant, perspective control. But you can also set its parameters to 768×768 or 512×512. Here are the best prompts for Stable Diffusion XL collected from the community on Reddit and Discord: 📷 Feb 24, 2024 · Here’s my list of the best SDXL prompts. Also, for all the prompts below, I’ve purely used the SDXL 1. Learned from Midjourney, the manual tweaking is not needed, and users only need to focus on the prompts and images. May 5, 2024 · However, the effect of step count depends on the sampler chosen. 1. For example, if you’re asking for a picture of a happy dog, you should use a negative prompt, like “No sad dogs”. 1-768. The best prompts are detailed, specific, and well-structured to help the model realize your vision. Fashion Model. ai discord score_9, score_8_up, score_7_up, score_6_up, score_5_up, score_4_up, just describe what you want Aug 3, 2023 · The Stable Diffusion model SDXL 1. In Stable Diffusion, models determine what art style they can Nov 12, 2023 · Stable Diffusion Prompts for NFT Art (14 Prompt Examples) In this guide I'll show you how to use a generative AI model like Stable Diffusion XL to create NFT art and characters. Moving into detailed subject and scene description, the focus is on precision. This will let you run the model from your PC. Text prompts are used as the primary conditioning factor. Conversely, a word inside square brackets receives reduced prominence. Read the SDXL guide for a more detailed walkthrough of how to use this model, and other techniques it uses to produce high quality images. Stable UnCLIP 2. Stable Diffusion is similarly powerful to DALL-E 2, but open source, and open to the public through Dream Studio, where anyone gets 50 free uses just by signing up with an email address. It's a versatile model that can generate diverse Stable Diffusion XL (SDXL) SDXL is a more powerful version of the Stable Diffusion model. And with the built-in styles, it’s much easier to control the output. 0 prompts and negative prompts to create stunning images with different styles. Renowned for its versatility, this model is adept in many scenarios, offering impressive performance across diverse applications. Stable Diffusion XL works especially well with images between 768 and 1024. Insights on Prompting SDXL Embarking on the SDXL model journey requires an understanding of its intrinsic nuances. g. It has a base resolution of 1024x1024 pixels. Try out this prompt on OpenArt Model: Stable Diffusion Seed: 70455 Scale: 13 Steps: 25 Resolution: 512 x 512 Sampler: DDIM 8 Learn how to create and use prompts for Stable Diffusion, ChatGPT and Midjourney models with PromptHero, the ultimate prompt search engine. Dec 27, 2022 · Well, you need to specify that. However, a great prompt can go a long way in generating the best output. 3. 0 is Stable Diffusion's next-generation model. Conclusion . The UNext is 3x larger. FlashAttention: XFormers flash attention can optimize your model even further with more speed and memory improvements. If not defined, one has to pass prompt_embeds. It’s significantly better than previous Stable Diffusion models at realism. Similarly, with Invoke AI, you just select the new sdxl model. Concise: Use concise language and avoid unnecessary words that may confuse the model or dilute the intended meaning. Steps: 4-6. 1. Redmond – App Icons Lora for SD XL 1. Mar 19, 2024 · Tips for good prompts. prompt #1: character covered in a lush, living tapestry of flora, with vines and leaves that appear as a natural, enchanted armor. Initially, the base model is deployed to establish the overall composition of the image. 0” is a LoRA that is fine tuned to create app icons sized for iOS and Android apps. As a ballpark, most samplers should use around 20 to 40 steps for the best balance between quality and speed. In Stable Diffusion, wrapping a word with triple parentheses ( ( (word))) boosts its weight by 1. Stable Diffusion separates the imaging process into a diffusion process at runtime. This is where you'll find your prompt and negative prompt. GBJI • 1 min. Stable Diffusion XL (SDXL) is a larger and more powerful iteration of the Stable Diffusion model, capable of producing higher resolution images. 35 Denoise + 1. 0 model. It’s worth mentioning that previous Blog post about Stable Diffusion: In-detail blog post explaining Stable Diffusion. CFG 7 – 10: Recommended for most prompts. I wanted to change the style for two prompts. While having an overview is helpful, keep in mind that these styles only imitate certain aspects of the artist's work (color, medium, location, etc. Consider aspects such as the subject matter, setting, mood, color scheme, and lighting. 1 comment. CFG 10 – 15: When you’re sure that your prompt is detailed and very clear on what you want the image to look like. You need to know that the model is a switchable part of AI where magic is stored. Mar 8, 2024 · This guide elucidates prime prompt tips and embodies 15 unique SDXL prompts that span a broad stylistic spectrum. Prompt weighting basics. Mar 7, 2024 · python3 client. As we all know, StabilityAI claims that the model is optimized for generating images from concise prompts. You can find a detailed guide on integrating our API's here: Stable Diffusion XL 1. This visualization forms the foundation of your prompt. Moreover, I May 27, 2023 · En outre, les prompts négatifs sont particulièrement importants dans les modèles Stable Diffusion v2. a masterpiece of an ocean with waves and palm trees, sunset in In this report, we discuss the art and science of prompt engineering for diffusion models primarily focussing on the Stable Diffusion family of models. Here, the use of text weights in prompts becomes important, allowing for emphasis on certain elements within the scene. py --model stable_diffusion_xl --prompt "butterfly in new york, 4k, realistic" --clients 2 –requests 5 Examine the server log and metrics. Software. You can use the syntax (keyword:weight) to control the weight of the keyword. In every step, the U-net in Stable Diffusion will use the prompt to guide the refinement of noise into a picture. Note that I use the Clipdrop platform by Stability AI for access to Stable Diffusion XL. CLIP can be used to create detailed prompts for stable diffusion models by providing relevant tags that describe the content of an image. 1 and 2 at the end. So, it’s like giving a little The Stable Diffusion XL (SDXL) model effectively comprises two distinct models working in tandem: ‍. Apr 3, 2024 · Here in our prompt, I used “3D Rendering” as my medium. It might take a few minutes to load the model fully. The higher resolution enables far greater detail and clarity in generated imagery. Since it uses SD XL 1. They are limited by the rather superficial knowledge of SD, but can probably give you a Jan 30, 2024 · Stable Diffusion is a text-to-image model, powered by AI, that uses deep learning to generate high-quality images from text. It saves you time and is great for quickly fixing common issues like garbled faces. Deforum is a popular way to make a video from a text prompt. Prompts are important because they describe what you want a diffusion model to generate. Dec 24, 2023 · Stable Diffusion XL (SDXL) is a powerful text-to-image generation model. In painting, the posture of a character is crucial for vividly expressing their life situation, emotional state, and mood. In this tutorial we will be going through the prompt formula for Stable Diffusion XL TURBO model and different parameter settings that will help you generate Aug 25, 2023 · No longer must users resort to terms like “masterpiece” to yield high-quality results. CFG: 1-2 (recommend 2 for a bit negative prompt affection) Negative: Only working slightly on CFG 2. instead. The Web UI offers various features, including generating images from text prompts (txt2img), image-to-image processing (img2img Jan 31, 2024 · Stable Diffusion Illustration Prompts. Intense, contemplative look, blending modern fashion with classic art aesthetics. If you want to run Stable Diffusion locally, you can follow these simple steps. 0_0. SDXL’s UNet is 3x larger and the model adds a second text encoder to the architecture. Unlikely that your image will look like their image. Apocalyptic Space Station. Whether you're looking to visualize concepts, explore new creative avenues, or enhance The training script is also similar to the Text-to-image training guide, but it’s been modified to support SDXL training. Step 1. Essentially, when a word is enclosed in parentheses, the model emphasizes it more in its output. Stable Diffusion is not like Midjourney or other popular image generation software, you can't just ask it what you want. If you’re eager to dive into the world of AI-generated art using Stable Diffusion XL, this guide will help you get started. The CLIP Text Enode node first converts the prompt into tokens and then encodes them into embeddings with the text encoder. E. Use an appropriate medium type consistent with the artist. Option 2: Install the extension stable-diffusion-webui-state. These prompts guide the image generation process, ensuring that the produced images align with the text's description. It wasn't long before many celebrities and Sep 8, 2023 · In this article, you will find some prompt tips and 15 SDXL prompts with different styles. The artist’s name is a very strong style modifier. Par exemple, l'ajout de "ugly" (laid), "disformed” (difforme) comme negative prompt peut aider à obtenir des images plus Feb 20, 2024 · Stable Diffusion Prompts Examples. A separate Refiner model based on Latent has been Jul 22, 2023 · After Detailer (adetailer) is a Stable Diffusion Automatic11111 web-UI extension that automates inpainting and more. Dreambooth - Quickly customize the model by fine-tuning it. The default weight = 1. But if you need to discover more image styles, you can check out this list where I covered 80+ Stable Diffusion styles. 9vae. To begin, envision the image you wish to create. Actually, It helps the generator understand what to avoid while creating the image. Feb 11, 2024 · Folders and source model Source model: sd_xl_base_1. prompt: “📸 Portrait of an aged Asian warrior chief 🌟, tribal panther makeup 🐾, side profile, intense gaze 👀, 50mm portrait photography 📷, dramatic rim lighting 🌅 –beta –ar 2:3 –beta –upbeta –upbeta”. I generate images at 1024×1024 and resize them to speed up my webpage. 8: Look at other people's prompts. Here, for each type of photorealistic image, I will first provide a brief explanation of that type. General info on Stable Diffusion - Info on other tasks that are powered by Stable Sep 22, 2023 · Option 1: Every time you generate an image, this text block is generated below your image. 9 at the end. gv bk pv pq qm wk iq xt uv ug